MODEL PROFILE

animagine-xl-4.0

cagliostrolab · animagine-xl-4.0

**Animagine XL 4.0**, also stylized as **Anim4gine**, is the ultimate anime-themed finetuned SDXL model and the latest installment of Animagine XL series. Despite being a continuation, the model was retrained from Stable Diffusion XL 1.0 with a massive dataset of 8.4M diverse anime-style images from various sources with the knowledge cut-off of January 7th 2025 and finetuned for approximately 2650 GPU hours. Similar to the previous version, this model was trained using tag ordering method for the identity and style training. With the release of **Animagine XL 4.0 Opt (Optimized)**, the model has been further refined with an additional dataset, improving **stability**, **anatomy accuracy**, **noise reduction**, **color saturation**, and **overall color accuracy**. These enhancements make **Animagine XL 4.0 Opt** more consistent and visually appealing while maintaining the signature quality of the series. - 2025-02-13 – Added Animagine XL 4.0 Opt - Better stability for more consistent outputs - Enhanced anatomy with more accurate proportions - Reduced noise and artifacts in generations - Fixed low saturation issues, resulting in richer colors - Improved color accuracy for more visually appealing results - 2025-01-24 – Initial release - **Developed by**: Cagliostro Research Lab - **Model type**: Diffusion-based text-to-image generative model - **License**: CreativeML Open RAIL++-M - **Model Description**: This is a model that can be used to generate and modify specifically anime-themed images based on text prompt - **Fine-tuned from**: Stable Diffusion XL 1.0

Image generationlicense:openrail++
Context windowNot disclosed
Input → OutputText → Image
Model weightsOpen weights
Profile updated2026/11/2 19:16:26
OVERVIEW

Model overview

**Animagine XL 4.0**, also stylized as **Anim4gine**, is the ultimate anime-themed finetuned SDXL model and the latest installment of Animagine XL series. Despite being a continuation, the model was retrained from Stable Diffusion XL 1.0 with a massive dataset of 8.4M diverse anime-style images from various sources with the knowledge cut-off of January 7th 2025 and finetuned for approximately 2650 GPU hours. Similar to the previous version, this model was trained using tag ordering method for the identity and style training. With the release of **Animagine XL 4.0 Opt (Optimized)**, the model has been further refined with an additional dataset, improving **stability**, **anatomy accuracy**, **noise reduction**, **color saturation**, and **overall color accuracy**. These enhancements make **Animagine XL 4.0 Opt** more consistent and visually appealing while maintaining the signature quality of the series. - 2025-02-13 – Added Animagine XL 4.0 Opt - Better stability for more consistent outputs - Enhanced anatomy with more accurate proportions - Reduced noise and artifacts in generations - Fixed low saturation issues, resulting in richer colors - Improved color accuracy for more visually appealing results - 2025-01-24 – Initial release - **Developed by**: Cagliostro Research Lab - **Model type**: Diffusion-based text-to-image generative model - **License**: CreativeML Open RAIL++-M - **Model Description**: This is a model that can be used to generate and modify specifically anime-themed images based on text prompt - **Fine-tuned from**: Stable Diffusion XL 1.0

Key features

特点摘要待依据来源补齐,不凭模型名称推断能力。

Limitations & considerations

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PRICING & COST

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INDEPENDENT BENCHMARKS

Public benchmarks & run configuration

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Frequently asked questions

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Sources & traceability

许可证:openrail++。模型介绍与第三方评测均以原来源的具体版本和方法为准。